SDXL:改进潜在扩散模型以实现高分辨率图像合成

SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis

罗宾·罗姆巴赫 Robin Rombach · Stability AI · 2023-07-04 · arXiv:2307.01952 ↗ · 被引 4946

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摘要 · Abstract

我们提出了 SDXL,一种用于文本到图像合成的潜在扩散模型。与之前版本的 Stable Diffusion 相比,SDXL 利用了大了三倍的 UNet 骨干网络:模型参数的增加主要归因于更多的注意力块和更大的交叉注意力上下文,因为 SDXL 使用了第二个文本编码器。我们设计了多种新颖的条件化方案,并在多个宽高比上训练 SDXL。我们还引入了一个精炼模型,该模型通过事后图像到图像技术用于提高 SDXL 生成样本的视觉保真度。我们证明,与之前版本的 Stable Diffusion 相比,SDXL 显示出显著改进的性能,并达到了与黑盒最先进图像生成器竞争的结果。为了促进开放研究和提高大型模型训练与评估的透明度,我们在 https://github.com/Stability-AI/generative-models 上提供代码和模型权重。

We present SDXL, a latent diffusion model for text-to-image synthesis. Compared to previous versions of Stable Diffusion, SDXL leverages a three times larger UNet backbone: The increase of model parameters is mainly due to more attention blocks and a larger cross-attention context as SDXL uses a second text encoder. We design multiple novel conditioning schemes and train SDXL on multiple aspect ratios. We also introduce a refinement model which is used to improve the visual fidelity of samples generated by SDXL using a post-hoc image-to-image technique. We demonstrate that SDXL shows drastically improved performance compared the previous versions of Stable Diffusion and achieves results competitive with those of black-box state-of-the-art image generators. In the spirit of promoting open research and fostering transparency in large model training and evaluation, we provide access to code and model weights at https://github.com/Stability-AI/generative-models

核心贡献 · Key contributions

局限 · Limitations

论文章节 · Sections(共 19)

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